Search Results for "diffusion"
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Showing 10 open source projects for "diffusion"
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Dream Textures
Stable Diffusion built-in to Blender
...Inpaint to fix up images and convert existing textures into seamless ones automatically. Outpaint to increase the size of an image by extending it in any direction. Perform style transfer and create novel animations with Stable Diffusion as a post processing step. Dream Textures has been tested with CUDA and Apple Silicon GPUs. Over 4GB of VRAM is recommended. -
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PersonaLive
Expressive Portrait Image Animation for Live Streaming
PersonaLive is an open-source diffusion-based portrait animation framework focused on generating expressive, long-duration animated sequences in real time, primarily for live streaming or interactive applications. It leverages deep generative models that condition on a static reference image and a driving input (such as motion or expression cues) to produce a seamless animated portrait sequence that can run indefinitely without segmentation artifacts. -
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Lama Cleaner
Image inpainting tool powered by SOTA AI Model
Image inpainting tool powered by SOTA AI Model. Remove any unwanted object, defect, or people from your pictures or erase and replace(powered by stable diffusion) anything on your pictures. Lama Cleaner is a free, open-source and fully self-hostable inpainting tool powered by state-of-the-art AI models. You can use it to remove any unwanted object, defect, or people from your pictures or erase and replace anything on your pictures. Many AICG creators are using Lama Cleaner to clean-up their work. ... -
4
IOPaint
Image inpainting tool powered by SOTA AI Model
IOPaint is a powerful open-source image editing tool focused on inpainting, outpainting, object removal, and general image manipulation driven by state-of-the-art AI models, delivering these capabilities through both local and hosted workflows. Designed to be fully self-hosted and flexible, IOPaint supports a variety of underlying generators and inpaint models — from LaMa erase networks to Stable Diffusion-based replace/object generation — giving users multiple ways to refine or reconstruct images by removing unwanted elements or expanding artwork beyond its original boundaries. Its feature set includes erasing people, watermarks, or defects, adding or replacing objects, applying text-aware edits, and extending images outward (outpainting) to fill contours or expand compositions. -
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5
Clarity AI Upscaler
AI Image Upscaler & Enhancer
Clarity AI Upscaler is an open-source AI image enhancement tool designed to increase the resolution and visual quality of images using modern generative techniques. The system uses deep learning models based on diffusion and other image generation methods to reconstruct high-resolution versions of low-resolution images while preserving important visual details. Unlike traditional interpolation-based upscaling algorithms, the system generates additional visual information that improves perceived clarity and sharpness. The project is intended as a free and open alternative to commercial AI upscaling tools, allowing developers and digital artists to run the technology locally or integrate it into their own workflows. ... -
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AnimateDiff
Plug-n-play module turning text-to-image models into animation
AnimateDiff is an open-source project designed to enhance text-to-image diffusion models by adding animation capabilities. It allows users to turn static images generated by popular text-to-image models into animated sequences without requiring additional model training. This plug-and-play tool is compatible with a wide range of community models and facilitates the generation of animation directly from pre-existing text-to-image models. -
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Stable Diffusion Rembg
Removes backgrounds from pictures. Extension for webui
This project is an extension for the Stable Diffusion Web UI that removes backgrounds from images directly inside the interface. It wraps popular background-removal models so creators can take a generated or uploaded image and isolate the subject with a single click. The workflow is designed to be non-destructive: you can preview, tweak thresholds, and export either a transparent PNG or a masked layer for further editing. -
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Stable Diffusion in Docker
Run the Stable Diffusion releases in a Docker container
Run the Stable Diffusion releases in a Docker container with txt2img, img2img, depth2img, pix2pix, upscale4x, and inpaint. Run the Stable Diffusion releases on Huggingface in a GPU-accelerated Docker container. By default, the pipeline uses the full model and weights which requires a CUDA capable GPU with 8GB+ of VRAM. It should take a few seconds to create one image. -
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DALL-E 2 - Pytorch
Implementation of DALL-E 2, OpenAI's updated text-to-image synthesis
Implementation of DALL-E 2, OpenAI's updated text-to-image synthesis neural network, in Pytorch. The main novelty seems to be an extra layer of indirection with the prior network (whether it is an autoregressive transformer or a diffusion network), which predicts an image embedding based on the text embedding from CLIP. Specifically, this repository will only build out the diffusion prior network, as it is the best performing variant (but which incidentally involves a causal transformer as the denoising network) To train DALLE-2 is a 3 step process, with the training of CLIP being the most important. ... -
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10
audio-diffusion-pytorch
Audio generation using diffusion models, in PyTorch
A fully featured audio diffusion library, for PyTorch. Includes models for unconditional audio generation, text-conditional audio generation, diffusion autoencoding, upsampling, and vocoding. The provided models are waveform-based, however, the U-Net (built using a-unet), DiffusionModel, diffusion method, and diffusion samplers are both generic to any dimension and highly customizable to work on other formats.
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